Hugging face stable diffusion Optimum Optimum provides a Stable Diffusion pipeline compatible with both OpenVINO and ONNX Runtime . Stable Diffusion 3. See full list on github. . Discover amazing ML apps made by the community Spaces Jun 12, 2024 · Stable Diffusion 3 Medium is a Multimodal Diffusion Transformer that can generate images based on text prompts. Stable Diffusion 3 Medium Model Stable Diffusion 3 Medium is a Multimodal Diffusion Transformer (MMDiT) text-to-image model that features greatly improved performance in image quality, typography, complex prompt understanding, and resource-efficiency. ai/license. 8k. Model Details Model Description (SVD) Image-to-Video is a latent diffusion model trained to generate short video clips Stable Diffusion v2-1 Model Card This model card focuses on the model associated with the Stable Diffusion v2-1 model, codebase available here. 1-768. Batch: 32 x 8 x 2 x 4 = 2048 Introduction to Stable Diffusion. Discover amazing ML apps made by the community Spaces. stable-diffusion. First 595k steps regular training, then 440k steps of inpainting training at resolution 512x512 on “laion-aesthetics v2 5+” and 10% dropping of the text-conditioning to improve classifier-free classifier-free guidance sampling . This stable-diffusion-2 model is resumed from stable-diffusion-2-base (512-base-ema. ckpt) and trained for 150k steps using a v-objective on the same dataset. This chapter introduces the building blocks of Stable Diffusion which is a generative artificial intelligence (generative AI) model that produces unique photorealistic images from text and image prompts. More details on model performance across various devices, can be found here. Stable Diffusion is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input. This model card gives an overview of all available model checkpoints. Aug 22, 2022 · Stable Diffusion is a text-to-image latent diffusion model created by the researchers and engineers from CompVis, Stability AI and LAION. Please note: For commercial use, please refer to https://stability. This repository provides scripts to run Stable-Diffusion on Qualcomm® devices. 5 Large is a new version of the diffusion model for image generation, with improved stability and quality. 🤗 Diffusers is the go-to library for state-of-the-art pretrained diffusion models for generating images, audio, and even 3D structures of molecules. Download the weights sd-v1-4. App Files Files Community 20280 Refreshing. 225,000 steps at resolution 512x512 on "laion-aesthetics v2 5+" and 10 % dropping of the text-conditioning to improve classifier-free guidance sampling. ai Aug 22, 2022 · We've gone from the basic use of Stable Diffusion using 🤗 Hugging Face Diffusers to more advanced uses of the library, and we tried to introduce all the pieces in a modern diffusion system. Optimizer: AdamW. Oct 29, 2024 · Stable Diffusion 3. com Stable Diffusion 3. It is trained on 512x512 images from a subset of the LAION-5B database. Finetuning a diffusion model on new data and adding guidance. Stable Diffusion v1-5 Model Card Stable Diffusion is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input. 98. 5. See examples of image generation from text prompts and how to customize the pipeline parameters. We recommend to explore different hyperparameters to get the best results on your dataset. 5 Large is a Multimodal Diffusion Transformer (MMDiT) text-to-image model that features improved performance in image quality, typography, complex prompt understanding, and resource-efficiency. stable-diffusion-v1-2: Resumed from stable-diffusion-v1-1. Features Detailed feature showcase with images: Original txt2img and img2img modes; One click install and run script (but you still must install python and git) Outpainting; Inpainting; Color Sketch; Prompt Matrix; Stable Diffusion Upscale The Stable-Diffusion-v-1-4 checkpoint was initialized with the weights of the Stable-Diffusion-v-1-2 checkpoint and subsequently fine-tuned on 225k steps at resolution 512x512 on "laion-aesthetics v2 5+" and 10% dropping of the text-conditioning to improve classifier-free guidance sampling. Resumed for another 140k steps on 768x768 images. For more in-detail model cards, please have a look at the model repositories listed under Model Access . ckpt; sd-v1-4-full-ema. 1), and then fine-tuned for another 155k extra steps with punsafe=0. Stable UnCLIP 2. Latent diffusion applies the diffusion process over a lower dimensional latent space to reduce memory and compute complexity. Please note: This model is released under the Stability Community License. 0, and an estimated watermark probability < 0. 5-medium-gguf This is a model from the MagicPrompt series of models, which are GPT-2 models intended to generate prompt texts for imaging AIs, in this case: Stable Diffusion. Running on CPU Upgrade. ckpt Oct 30, 2023 · city96/stable-diffusion-3. art". 5 Medium is a Multimodal Diffusion Transformer with improvements (MMDiT-X) text-to-image model that features improved performance in image quality, typography, complex prompt understanding, and resource-efficiency. Learn how to use Stable Diffusion, a text-to-image latent diffusion model, with the Diffusers library. Follow the steps to create an endpoint, test and generate images, and integrate the model via API with Python. Learn how to use it with Diffusers, a library for working with Hugging Face's models and pipelines. March 24, 2023. 5-large-turbo-gguf. ckpt) with an additional 55k steps on the same dataset (with punsafe=0. For more information about how Stable Diffusion functions, please have a look at 🤗's Stable Diffusion blog. Jun 12, 2024 · Stable Diffusion 3 Medium is a Multimodal Diffusion Transformer (MMDiT) text-to-image model that features greatly improved performance in image quality, typography, complex prompt understanding, and resource-efficiency. like 10. It’s easy to overfit and run into issues like catastrophic forgetting. and get access to the augmented documentation experience Please visit this very in-detail blog post on Stable Diffusion! This model is an implementation of Stable-Diffusion found here. If you liked this topic and want to learn more, we recommend the following resources: We’re on a journey to advance and democratize artificial intelligence through open source and open science. Latent diffusion applies the diffusion process over a lower dimensional latent space to reduce memory and compute complexity. 🖼️ Here's an example: This model was trained with 150,000 steps and a set of about 80,000 data filtered and extracted from the image finder for Stable Diffusion: "Lexica. Stable Diffusion web UI A browser interface based on Gradio library for Stable Diffusion. It is a free research model for non-commercial and commercial use, with different variants and text encoders available. Stable Video Diffusion (SVD) Image-to-Video is a diffusion model that takes in a still image as a conditioning frame, and generates a video from it. Image-to-image. Stable Diffusion is a text-to-image latent diffusion model created by the researchers and engineers from CompVis, Stability AI and LAION. Text-to-Image • Updated Oct 23 • 4. For some workflow examples and see what ComfyUI can do you can check out: ComfyUI Examples Installing ComfyUI Features For more information on how to use Stable Diffusion XL with diffusers, please have a look at the Stable Diffusion XL Docs. The Stable-Diffusion-Inpainting was initialized with the weights of the Stable-Diffusion-v-1-2. 19 A powerful and modular stable diffusion GUI and backend. New stable diffusion finetune (Stable unCLIP 2. This stable-diffusion-2-1 model is fine-tuned from stable-diffusion-2 (768-v-ema. Nov 28, 2022 · Learn how to deploy and use Stable Diffusion, a text-to-image latent diffusion model, on Hugging Face Inference Endpoints. 1, Hugging Face) at 768x768 resolution, based on SD2. stabilityai / stable-diffusion. stable-diffusion-v1-4 Resumed from stable-diffusion-v1-2. The Stable Diffusion model can also be applied to image-to-image generation by passing a text prompt and an initial image to condition the generation of new images. Join the Hugging Face community. The text-to-image fine-tuning script is experimental. App Files Files Community 20280 Since its public release the community has done an incredible job at working together to make the stable diffusion checkpoints faster, more memory efficient, and more performant. Gradient Accumulations: 2. Model Access Each checkpoint can be used both with Hugging Face's 🧨 Diffusers library or the original Stable Diffusion GitHub repository. Stable Diffusion pipelines. 🧨 Diffusers offers a simple API to run stable diffusion with all memory, computing, and quality improvements. This ui will let you design and execute advanced stable diffusion pipelines using a graph/nodes/flowchart based interface. This model allows for image variations and mixing operations as described in Hierarchical Text-Conditional Image Generation with CLIP Latents, and, thanks to its modularity, can be combined with other models such as KARLO. For more technical details, please refer to the Research paper. 515,000 steps at resolution 512x512 on "laion-improved-aesthetics" (a subset of laion2B-en, filtered to images with an original size >= 512x512, estimated aesthetics score > 5. stable-diffusion. Unit 3: Stable Diffusion Exploring a powerful text-conditioned latent diffusion model; Unit 4: Doing more with diffusion Advanced techniques for going further with diffusion; Who are we? About the authors: Jonathan Whitaker is a Data Scientist/AI Researcher doing R&D with answer. Hardware: 32 x 8 x A100 GPUs. 5 Medium Model Stable Diffusion 3. 15k • 35 city96/stable-diffusion-3. 1. Model Details Model Type: Image generation; Model Stats: Input: Text prompt to generate image; QNN-SDK: 2. ifdu csoh oxbzzd uvly owcb netf fwadj peqlvqt delz itms